Cocktail is a standalone desktop app that uses the Civitai API combined with a local database to. 1 to make it work you need to use . More up to date and experimental versions available at: Results oversaturated, smooth, lacking detail? No. articles. 75, Hires upscale: 2, Hires steps: 40, Hires upscaler: Latent (bicubic antialiased) Most of the sample images are generated with hires. Install Stable Diffusion Webui's Extension tab, go to Install from url sub-tab. Epîc Diffusion is a general purpose model based on Stable Diffusion 1. Try to experiment with the CFG scale, 10 can create some amazing results but to each their own. This model has been trained on 26,949 high resolution and quality Sci-Fi themed images for 2 Epochs. 2. Highres-fix (upscaler) is strongly recommended (using the SwinIR_4x,R-ESRGAN 4x+anime6B by. mutsuki_mix. Trigger is arcane style but I noticed this often works even without it. This model is a 3D merge model. The official SD extension for civitai takes months for developing and still has no good output. . Plans Paid; Platforms Social Links Visit Website Add To Favourites. These models perform quite well in most cases, but please note that they are not 100%. x intended to replace the official SD releases as your default model. This one's goal is to produce a more "realistic" look in the backgrounds and people. Do check him out and leave him a like. 0. Fix. Through this process, I hope not only to gain a deeper. Which includes characters, background, and some objects. and, change about may be subtle and not drastic enough. Facbook Twitter linkedin Copy link. 5 (general), 0. He was already in there, but I never got good results. Once you have Stable Diffusion, you can download my model from this page and load it on your device. Unlike other anime models that tend to have muted or dark colors, Mistoon_Ruby uses bright and vibrant colors to make the characters stand out. Activation words are princess zelda and game titles (no underscores), which I'm not gonna list, as you can see them from the example prompts. Cherry Picker XL. Sampler: DPM++ 2M SDE Karras. . Face restoration is still recommended. Join our 404 Contest and create images to populate our 404 pages! Running NOW until Nov 24th. 2 and Stable Diffusion 1. Support☕ more info. Hope you like it! Example Prompt: <lora:ldmarble-22:0. It is advisable to use additional prompts and negative prompts. Created by ogkalu, originally uploaded to huggingface. Stable Diffusion Models, sometimes called checkpoint models, are pre-trained Stable Diffusion weights for generating a particular style of images. You can use some trigger words (see Appendix A) to generate specific styles of images. . 2发布，用DARKTANG融合REALISTICV3版Human Realistic - Realistic V. For v12_anime/v4. merging another model with this one is the easiest way to get a consistent character with each view. 8, but weights from 0. Look at all the tools we have now from TIs to LoRA, from ControlNet to Latent Couple. Use the negative prompt: "grid" to improve some maps, or use the gridless version. Guaranteed NSFW or your money back Fine-tuned from Stable Diffusion v2-1-base 19 epochs of 450,000 images each, co. 0 updated. SDXLベースモデルなので、SD1. SafeTensor. 1 to make it work you need to use . Just enter your text prompt, and see the generated image. 5 as well) on Civitai. Used to named indigo male_doragoon_mix v12/4. Pony Diffusion is a Stable Diffusion model that has been fine-tuned on high-quality pony, furry and other non photorealistic SFW and NSFW images. Different models available, check the blue tabs above the images up top: Stable Diffusion 1. Refined v11. 05 23526-1655-下午好. The recommended sampling is k_Euler_a or DPM++ 2M Karras on 20 steps, CFGS 7. This checkpoint includes a config file, download and place it along side the checkpoint. Action body poses. 5 Beta 3 is fine-tuned directly from stable-diffusion-2-1 (768), using v-prediction and variable aspect bucketing (maximum pixel. Once you have Stable Diffusion, you can download my model from this page and load it on your device. v8 is trash. The Stable Diffusion 2. It's a more forgiving and easier to prompt SD1. I don't remember all the merges I made to create this model. . LORA: For anime character LORA, the ideal weight is 1. If using the AUTOMATIC1111 WebUI, then you will. 0 significantly improves the realism of faces and also greatly increases the good image rate. Resource - Update. Except for one. You will need the credential after you start AUTOMATIC11111. posts. A startup called Civitai — a play on the word Civitas, meaning community — has created a platform where members can post their own Stable Diffusion-based AI. Likewise, it can work with a large number of other lora, just be careful with the combination weights. Another LoRA that came from a user request. It tends to lean a bit towards BoTW, but it's very flexible and allows for most Zelda versions. 8 is often recommended. If you are the person or a legal representative of the person depicted, and would like to request the removal of this resource, you can do so here. A lot of checkpoints available now are mostly based on anime illustrations oriented towards 2. V1 (main) and V1. You can view the final results with. Shinkai Diffusion is a LORA trained on stills from Makoto Shinkai's beautiful anime films made at CoMix Wave Films. This model has been archived and is not available for download. Saves on vram usage and possible NaN errors. Install the Civitai Extension: Begin by installing the Civitai extension for the Automatic 1111 Stable Diffusion Web UI. If you are the person or a legal representative of the person depicted, and would like to request the removal of this resource, you can do so here. • 9 mo. . To utilize it, you must include the keyword " syberart " at the beginning of your prompt. Vampire Style. merging another model with this one is the easiest way to get a consistent character with each view. In the second step, we use a. This is a fine-tuned Stable Diffusion model (based on v1. In my tests at 512,768 resolution, the good image rate of the Prompts I used before was above 50%. V1: A total of ~100 training images of tungsten photographs taken with CineStill 800T were used. Browse nsfw Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAsRecommend: DPM++2M Karras, Clip skip 2 Sampler, Steps: 25-35+. 5 for generating vampire portraits! Using a variety of sources such as movies, novels, video games, and cosplay photos, I've trained the model to produce images with all the classic vampire features like fangs and glowing eyes. Give your model a name and then select ADD DIFFERENCE (This will make sure to add only the parts of the inpainting model that will be required) Select ckpt or safetensors. Essentials extensions and settings for Stable Diffusion for the use with Civit AI. Classic NSFW diffusion model. Use it at around 0. 0 (B1) Status (Updated: Nov 18, 2023): - Training Images: +2620 - Training Steps: +524k - Approximate percentage of completion: ~65%. For next models, those values could change. Cmdr2's Stable Diffusion UI v2. KayWaii. Clip Skip: It was trained on 2, so use 2. このモデルは3D系のマージモデルです。. 2版本时，可以. 9). Original Hugging Face Repository Simply uploaded by me, all credit goes to . 5 and 10 CFG Scale and between 25 and 30 Steps with DPM++ SDE Karras. If your characters are always wearing jackets/half off jackets, try adding "off shoulder" in negative prompt. Simple LoRA to help with adjusting a subjects traditional gender appearance. SDXL consists of a two-step pipeline for latent diffusion: First, we use a base model to generate latents of the desired output size. If you don't like the color saturation you can decrease it by entering oversaturated in negative prompt. Stable Diffusion Webui Extension for Civitai, to help you handle models much more easily. V3. The resolution should stay at 512 this time, which is normal for Stable Diffusion. Analog Diffusion. stable-diffusion. Join us on our Discord: collection of OpenPose skeletons for use with ControlNet and Stable Diffusion. 0 significantly improves the realism of faces and also greatly increases the good image rate. Style model for Stable Diffusion. Civitai Helper 2 also has status news, check github for more. 6-0. . Navigate to Civitai: Open your web browser, type in the Civitai website’s address, and immerse yourself. So far so good for me. 1 and Exp 7/8, so it has its unique style with a preference for Big Lips (and who knows what else, you tell me). BeenYou - R13 | Stable Diffusion Checkpoint | Civitai. The GhostMix-V2. . Prompt suggestions :use cartoon in prompt for more cartoonish images, you can use anime or realistic prompts both works the same. When applied, the picture will look like the character is bordered. 4 - a true general purpose model, producing great portraits and landscapes. Updated - SECO: SECO = Second-stage Engine Cutoff (I watch too many SpaceX launches!!) - am cutting this model off now, and there may be an ICBINP XL release, but will see what happens. The Civitai Discord server is described as a lively community of AI art enthusiasts and creators. Sensitive Content. Copy the file 4x-UltraSharp. Originally posted by nousr on HuggingFaceOriginal Model Dpepteahand3. This model was trained on the loading screens, gta storymode, and gta online DLCs artworks. These files are Custom Workflows for ComfyUI. Kenshi is my merge which were created by combining different models. 5d, which retains the overall anime style while being better than the previous versions on the limbs, but the light and shadow and lines are more like 2. Civitai stands as the singular model-sharing hub within the AI art generation community. I had to manually crop some of them. 3 | Stable Diffusion Checkpoint | Civitai，相比前作REALTANG刷图评测数据更好testing (civitai. I apologize as the preview images for both contain images generated with both, but they do produce similar results, try both and see which works. AS-Elderly: Place at the beginning of your positive prompt at strength of 1. Western Comic book styles are almost non existent on Stable Diffusion. HERE! Photopea is essentially Photoshop in a browser. And full tutorial on my Patreon, updated frequently. The model files are all pickle-scanned for safety, much like they are on Hugging Face. It is a Latent Diffusion Model that uses two fixed, pretrained text encoders ( OpenCLIP-ViT/G and CLIP-ViT/L ). (safetensors are recommended) And hit Merge. pth. stable Diffusion models, embeddings, LoRAs and more. I'm just collecting these. Since this embedding cannot drastically change the artstyle and composition of the image, not one hundred percent of any faulty anatomy can be improved. pit next to them. Space (main sponsor) and Smugo. C:stable-diffusion-uimodelsstable-diffusion)Redshift Diffusion. Trained on 70 images. sassydodo. fix is needed for prompts where the character is far away in order to make decent images, it drastically improve the quality of face and eyes! Sampler: DPM++ SDE Karras: 20 to 30 steps. The first version I'm uploading is a fp16-pruned with no baked vae, which is less than 2 GB, meaning you can get up to 6 epochs in the same batch on a colab. Review username and password. The purpose of DreamShaper has always been to make "a better Stable Diffusion", a model capable of doing everything on its own, to weave dreams. 1. These are optional files, producing similar results to the official ControlNet models, but with added Style and Color functions. Please consider to support me via Ko-fi. This upscaler is not mine, all the credit go to: Kim2091 Official WiKi Upscaler page: Here License of use it: Here HOW TO INSTALL: Rename the file from: 4x-UltraSharp. 5D/3D images) Steps : 30+ (I strongly suggest 50 for complex prompt) AnimeIllustDiffusion is a pre-trained, non-commercial and multi-styled anime illustration model. It's a model that was merged using a supermerger ↓↓↓ fantasticmix2. In addition, although the weights and configs are identical, the hashes of the files are different. Enable Quantization in K samplers. This checkpoint recommends a VAE, download and place it in the VAE folder. These first images are my results after merging this model with another model trained on my wife. Should work well around 8-10 cfg scale and I suggest you don't use the SDXL refiner, but instead do a i2i step on the upscaled. This is a fine-tuned variant derived from Animix, trained with selected beautiful anime images. The third example used my other lora 20D. Each pose has been captured from 25 different angles, giving you a wide range of options. Due to plenty of contents, AID needs a lot of negative prompts to work properly. Copy image prompt and setting in a format that can be read by Prompts from file or textbox. This is the first model I have published, and previous models were only produced for internal team and partner commercial use. All models, including Realistic Vision. Download (2. Overview. prompts that i always add: award winning photography, Bokeh, Depth of Field, HDR, bloom, Chromatic Aberration ,Photorealistic,extremely detailed, trending on artstation, trending. " (mostly for v1 examples)Browse pixel art Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAs75T: The most ”easy to use“ embedding, which is trained from its accurate dataset created in a special way with almost no side effects. Paste it into the textbox below the webui script "Prompts from file or textbox". Waifu Diffusion - Beta 03. Simply copy paste to the same folder as selected model file. Fine-tuned LoRA to improve the effects of generating characters with complex body limbs and backgrounds. This is the fine-tuned Stable Diffusion model trained on images from the TV Show Arcane. This extension allows you to seamlessly. Conceptually middle-aged adult 40s to 60s, may vary by model, lora, or prompts. Sci-Fi Diffusion v1. The idea behind Mistoon_Anime is to achieve the modern anime style while keeping it as colorful as possible. . 75T: The most ”easy to use“ embedding, which is trained from its accurate dataset created in a special way with almost no side effects. Settings are moved to setting tab->civitai helper section. ℹ️ The Babes Kissable Lips model is based on a brand new training, that is mixed with Babes 1. Originally shared on GitHub by guoyww Learn about how to run this model to create animated images on GitHub. Space (main sponsor) and Smugo. I am a huge fan of open source - you can use it however you like with only restrictions for selling my models. It shouldn't be necessary to lower the weight. List of models. Recommend. Instead, use the "Tiled Diffusion" mode to enlarge the generated image and achieve a more realistic skin texture. The Civitai model information, which used to fetch real-time information from the Civitai site, has been removed. Dreamlike Diffusion 1. Hello my friends, are you ready for one last ride with Stable Diffusion 1. 1_realistic: Hello everyone! These two are merge models of a number of other furry/non furry models, they also have mixed in a lot. When comparing stable-diffusion-howto and civitai you can also consider the following projects: stable-diffusion-webui-colab - stable diffusion webui colab. It fits greatly for architectures. Instead, the shortcut information registered during Stable Diffusion startup will be updated. 5 model to create isometric cities, venues, etc more precisely. Thanks for using Analog Madness, if you like my models, please buy me a coffee ️ [v6. work with Chilloutmix, can generate natural, cute, girls. This model is available on Mage. For example, “a tropical beach with palm trees”. Upload 3. These first images are my results after merging this model with another model trained on my wife. 4 (unpublished): MothMix 1. It proudly offers a platform that is both free of charge and open. Some Stable Diffusion models have difficulty generating younger people. 1 to make it work you need to use . 5 using +124000 images, 12400 steps, 4 epochs +32 training hours. The purpose of DreamShaper has always been to make "a. I spent six months figuring out how to train a model to give me consistent character sheets to break apart in Photoshop and animate. Steps and CFG: It is recommended to use Steps from “20-40” and CFG scale from “6-9”, the ideal is: steps 30, CFG 8. Pixar Style Model. This model uses the core of the Defacta 3rd series, but has been largely converted to a realistic model. flip_aug is a trick to learn more evenly, as if you had more images, but makes the AI confuse left and right, so it's your choice. Trigger word: 2d dnd battlemap. Raising from the ashes of ArtDiffusionXL-alpha, this is the first anime oriented model I make for the XL architecture. Then go to your WebUI, Settings -> Stable Diffusion on the left list -> SD VAE, choose your downloaded VAE. 6-1. It has been trained using Stable Diffusion 2. yaml). Upscaler: 4x-Ultrasharp or 4X NMKD Superscale. 8346 models. Triggers with ghibli style and, as you can see, it should work. This embedding can be used to create images with a "digital art" or "digital painting" style. For instance: On certain image-sharing sites, many anime character LORAs are overfitted. Other upscalers like Lanczos or Anime6B tends to smoothen them out, removing the pastel-like brushwork. Whether you are a beginner or an experienced user looking to study the classics, you are in the right place. Resources for more information: GitHub. This checkpoint includes a config file, download and place it along side the checkpoint. If you want to suppress the influence on the composition, please. py file into your scripts directory. Soda Mix. So veryBadImageNegative is the dedicated negative embedding of viewer-mix_v1. Avoid anythingv3 vae as it makes everything grey. huggingface. NAI is a model created by the company NovelAI modifying the Stable Diffusion architecture and training method. yaml). In simple terms, inpainting is an image editing process that involves masking a select area and then having Stable Diffusion redraw the area based on user input. . 5 as w. Asari Diffusion. Add a ️ to receive future updates. V7 is here. Most of the sample images follow this format. Browse weapons Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAsA dreambooth-method finetune of stable diffusion that will output cool looking robots when prompted. . Research Model - How to Build Protogen ProtoGen_X3. Step 2. Originally posted to HuggingFace by leftyfeep and shared on Reddit. Review Save_In_Google_Drive option. Stable Diffusion Webui Extension for Civitai, to help you handle models much more easily. NOTE: usage of this model implies accpetance of stable diffusion's CreativeML Open. Huggingface is another good source though the interface is not designed for Stable Diffusion models. nudity) if. Submit your Part 2 Fusion images here, for a chance to win $5,000 in prizes!Created by Astroboy, originally uploaded to HuggingFace. Now the world has changed and I’ve missed it all. This is a fine-tuned variant derived from Animix, trained with selected beautiful anime images. Speeds up workflow if that's the VAE you're going to use anyway. 5 version please pick version 1,2,3 I don't know a good prompt for this model, feel free to experiment i also have. I recommend you use an weight of 0. g. Counterfeit-V3 (which has 2. >Adetailer enabled using either 'face_yolov8n' or. All the images in the set are in png format with the background removed, making it possible to use multiple images in a single scene. ℹ️ The core of this model is different from Babes 1. This model performs best in the 16:9 aspect ratio, although it can also produce good results in a square format. 6. Posted first on HuggingFace. 5D/3D images) Steps : 30+ (I strongly suggest 50 for complex prompt)AnimeIllustDiffusion is a pre-trained, non-commercial and multi-styled anime illustration model. You must include a link to the model card and clearly state the full model name (Perpetual Diffusion 1. >Initial dimensions 512x615 (WxH) >Hi-res fix by 1. This might take some time. Civitai Helper 2 also has status news, check github for more. The last sample image shows a comparison between three of my mix models: Aniflatmix, Animix, and Ambientmix (this model). Warning: This model is NSFW. This checkpoint includes a config file, download and place it along side the checkpoint. AingDiffusion (read: Ah-eeng Diffusion) is a merge of a bunch of anime models. Enter our Style Capture & Fusion Contest! Part 2 of our Style Capture & Fusion contest is running until November 10th at 23:59 PST. . 2. This model imitates the style of Pixar cartoons. Sticker-art. A fine tuned diffusion model that attempts to imitate the style of late '80s early 90's anime specifically, the Ranma 1/2 anime. To utilize it, you must include the keyword " syberart " at the beginning of your prompt. Stable Diffusion on syväoppimiseen perustuva tekoälyohjelmisto, joka tuottaa kuvia tekstimuotoisesta kuvauksesta. This checkpoint includes a config file, download and place it along side the checkpoint. ”. Use "80sanimestyle" in your prompt. Use silz style in your prompts. The official SD extension for civitai takes months for developing and still has no good output. It allows users to browse, share, and review custom AI art models, providing a space for creators to showcase their work and for users to find inspiration. yaml file with name of a model (vector-art. リアル系マージモデルです。 このマージモデルを公開するにあたり、使用したモデルの製作者の皆様に感謝申し上げます。 This is a realistic merge model. These are the concepts for the embeddings. . vae. Included 2 versions, 1 for 4500 steps which is generally good, and 1 with some added input images for ~8850 steps, which is a bit cooked but can sometimes provide results closer to what I was after. This is a general purpose model able to do pretty much anything decently well from realistic to anime to backgrounds All the images are raw outputs. You may need to use the words blur haze naked in your negative prompts. Realistic Vision V6. ranma_diffusion. For more information, see here . • 15 days ago. While we can improve fitting by adjusting weights, this can have additional undesirable effects. Please read this! How to remove strong. com/models/38511?modelVersionId=44457 的DDicon模型使用，生成玻璃质感web风格B端元素。 v1和v2版本建议对应使用，v1. No animals, objects or backgrounds. The comparison images are compressed to . Browse controlnet Stable Diffusion models, checkpoints, hypernetworks, textual inversions, embeddings, Aesthetic Gradients, and LORAsSeeing my name rise on the leaderboard at CivitAI is pretty motivating, well, it was motivating, right up until I made the mistake of running my mouth at the wrong mod, didn't realize that was a ToS breach, or that bans were even a thing,. Silhouette/Cricut style. Since I use A111. 增强图像的质量，削弱了风格。. I don't remember all the merges I made to create this model. 1 recipe, also it has been inspired a little bit by RPG v4. Clip Skip: It was trained on 2, so use 2. Just another good looking model with a sad feeling . Out of respect for this individual and in accordance with our Content Rules, only work-safe images and non-commercial use is permitted. Head to Civitai and filter the models page to “ Motion ” – or download from the direct links in the table above. articles. 1 (512px) to generate cinematic images. For some reasons, the model stills automatically include in some game footage, so landscapes tend to look. If you gen higher resolutions than this, it will tile the latent space. 5 weight. If you like the model, please leave a review! This model card focuses on Role Playing Game portrait similar to Baldur's Gate, Dungeon and Dragon, Icewindale, and more modern style of RPG character. Based on StableDiffusion 1. It merges multiple models based on SDXL. This resource is intended to reproduce the likeness of a real person. LORA: For anime character LORA, the ideal weight is 1. Example: knollingcase, isometic render, a single cherry blossom tree, isometric display case, knolling teardown, transparent data visualization infographic, high-resolution OLED GUI interface display, micro-details, octane render, photorealism, photorealistic. Black Area is the selected or "Masked Input". He is not affiliated with this. The GhostMix-V2. 7 here) >, Trigger Word is ' mix4 ' . 首先暗图效果比较好，dark合适. Prompts listed on left side of the grid, artist along the top. This is a checkpoint that's a 50% mix of AbyssOrangeMix2_hard and 50% Cocoa from Yohan Diffusion. 5 with Automatic1111's checkpoint merger tool (Couldn't remember exactly the merging ratio and the interpolation method)About This LoRA is intended to generate an undressed version of the subject (on the right) alongside a clothed version (on the left). 45 GB) Verified: 14 days ago. Installation: As it is model based on 2. It does portraits and landscapes extremely well, animals should work too. Conceptually elderly adult 70s +, may vary by model, lora, or prompts. Noosphere - v3 | Stable Diffusion Checkpoint | Civitai. I am pleased to tell you that I have added a new set of poses to the collection. Then go to your WebUI, Settings -> Stable Diffusion on the left list -> SD VAE, choose your downloaded VAE. The Civitai Link Key is a short 6 character token that you'll receive when setting up your Civitai Link instance (you can see it referenced here in this Civitai Link installation video). 1. Sit back and enjoy reading this article whose purpose is to cover the essential tools needed to achieve satisfaction during your Stable Diffusion experience. Test model created by PublicPrompts This version contains a lot of biases but it does create a lot of cool designs of various subject will be creat. This LoRA model was finetuned on an extremely diverse dataset of 360° equirectangular projections with 2104 captioned training images, using the Stable Diffusion v1-5 model. The yaml file is included here as well to download. Copy the file 4x-UltraSharp. This is a simple extension to add a Photopea tab to AUTOMATIC1111 Stable Diffusion WebUI. PEYEER - P1075963156. Refined_v10. Copy as single line prompt. For example, “a tropical beach with palm trees”. . Originally uploaded to HuggingFace by NitrosockeThe new version is an integration of 2. This method is mostly tested on landscape. The version is not about the newer the better.